NexaAI/sdxl-base-int4-cuda

Stable-Diffusion-XL-Base-1.0

How to run

Visit sdk.nexa.ai/model

Model Description

Stable Diffusion XL Base 1.0 (SDXL 1.0) is a foundation text-to-image model released by Stability AI.
It is the flagship successor to Stable Diffusion 2.1, designed for photorealism, artistic flexibility, and high-resolution generation.

SDXL 1.0 is a latent diffusion model trained on a broad dataset of images and captions. Compared to prior versions, it improves prompt alignment, visual coherence, and output quality, especially in complex scenes and detailed compositions.

Features

• High fidelity image generation: sharper details and improved realism.
• Flexible style range: from photorealistic renders to artistic illustration.
• Better prompt alignment: improved understanding of nuanced or multi-concept prompts.
• High resolution support: natively trained for 1024×1024 images.
• Compositional strength: more accurate handling of multiple subjects and fine object placement.

Use Cases

• Creative content generation (illustrations, art, concept design)
• Product mockups and marketing visuals
• Character and environment ideation
• Storyboarding and visual storytelling
• Research in generative imaging

Inputs and Outputs

Input:

• Text prompts (descriptions, concepts, artistic directions)
• Optional negative prompts to avoid undesired elements

Output:

• Generated image(s) matching the prompt
• Default resolution: 1024×1024 pixels

License

• Licensed under: CreativeML Open RAIL++-M License

References

• Model card: https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0